mirror of
https://github.com/Sygil-Dev/sygil-webui.git
synced 2024-12-15 14:31:44 +03:00
96 lines
4.4 KiB
Markdown
96 lines
4.4 KiB
Markdown
# Web based UI for Stable Diffusion by [sd-webui](https://github.com/sd-webui)
|
|
|
|
## [Visit sd-webui's Discord Server](https://discord.gg/gyXNe4NySY) [](https://discord.gg/gyXNe4NySY)
|
|
|
|
## Installation instructions for [Windows](https://sd-webui.github.io/stable-diffusion-webui/docs/1.installation.html), [Linux](https://sd-webui.github.io/stable-diffusion-webui/docs/1.linux-installation.html)
|
|
|
|
### Want to ask a question or request a feature?
|
|
|
|
Come to our [Discord Server](https://discord.gg/gyXNe4NySY) or use [Discussions](https://github.com/sd-webui/stable-diffusion-webui/discussions).
|
|
|
|
## Documentation
|
|
|
|
[Documentaion is located here](https://sd-webui.github.io/stable-diffusion-webui/)
|
|
|
|
## Want to contribute?
|
|
|
|
Check the [Contribution Guide](CONTRIBUTING.md)
|
|
|
|
[sd-webui](https://github.com/sd-webui) is
|
|
*  [hlky](https://github.com/hlky)
|
|
* [ZeroCool940711](https://github.com/ZeroCool940711)
|
|
* [codedealer](https://github.com/codedealer)
|
|
|
|
|
|
|
|
## Gradio
|
|
|
|
### Features
|
|
|
|
### Screenshots
|
|
|
|
## Streamlit
|
|
|
|
### Features
|
|
|
|
### Screenshots
|
|
|
|
|
|
|
|
|
|
--------------
|
|
|
|
*Stable Diffusion was made possible thanks to a collaboration with [Stability AI](https://stability.ai/) and [Runway](https://runwayml.com/) and builds upon our previous work:*
|
|
|
|
[**High-Resolution Image Synthesis with Latent Diffusion Models**](https://ommer-lab.com/research/latent-diffusion-models/)<br/>
|
|
[Robin Rombach](https://github.com/rromb)\*,
|
|
[Andreas Blattmann](https://github.com/ablattmann)\*,
|
|
[Dominik Lorenz](https://github.com/qp-qp)\,
|
|
[Patrick Esser](https://github.com/pesser),
|
|
[Björn Ommer](https://hci.iwr.uni-heidelberg.de/Staff/bommer)<br/>
|
|
|
|
**CVPR '22 Oral**
|
|
|
|
which is available on [GitHub](https://github.com/CompVis/latent-diffusion). PDF at [arXiv](https://arxiv.org/abs/2112.10752). Please also visit our [Project page](https://ommer-lab.com/research/latent-diffusion-models/).
|
|
|
|
|
|
[Stable Diffusion](#stable-diffusion-v1) is a latent text-to-image diffusion
|
|
model.
|
|
Thanks to a generous compute donation from [Stability AI](https://stability.ai/) and support from [LAION](https://laion.ai/), we were able to train a Latent Diffusion Model on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) database.
|
|
Similar to Google's [Imagen](https://arxiv.org/abs/2205.11487),
|
|
this model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts.
|
|
With its 860M UNet and 123M text encoder, the model is relatively lightweight and runs on a GPU with at least 10GB VRAM.
|
|
See [this section](#stable-diffusion-v1) below and the [model card](https://huggingface.co/CompVis/stable-diffusion).
|
|
|
|
Stable Diffusion v1 refers to a specific configuration of the model
|
|
architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet
|
|
and CLIP ViT-L/14 text encoder for the diffusion model. The model was pretrained on 256x256 images and
|
|
then finetuned on 512x512 images.
|
|
|
|
*Note: Stable Diffusion v1 is a general text-to-image diffusion model and therefore mirrors biases and (mis-)conceptions that are present
|
|
in its training data.
|
|
Details on the training procedure and data, as well as the intended use of the model can be found in the corresponding [model card](https://huggingface.co/CompVis/stable-diffusion).
|
|
|
|
- Our codebase for the diffusion models builds heavily on [OpenAI's ADM codebase](https://github.com/openai/guided-diffusion)
|
|
and [https://github.com/lucidrains/denoising-diffusion-pytorch](https://github.com/lucidrains/denoising-diffusion-pytorch).
|
|
Thanks for open-sourcing!
|
|
|
|
- The implementation of the transformer encoder is from [x-transformers](https://github.com/lucidrains/x-transformers) by [lucidrains](https://github.com/lucidrains?tab=repositories).
|
|
|
|
|
|
BibTeX
|
|
|
|
```
|
|
@misc{rombach2021highresolution,
|
|
title={High-Resolution Image Synthesis with Latent Diffusion Models},
|
|
author={Robin Rombach and Andreas Blattmann and Dominik Lorenz and Patrick Esser and Björn Ommer},
|
|
year={2021},
|
|
eprint={2112.10752},
|
|
archivePrefix={arXiv},
|
|
primaryClass={cs.CV}
|
|
}
|
|
|
|
```
|
|
|
|
|